stable-diffusion-v-1-4-original
Maintainer: CompVis
2.7K
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| Property | Value |
|---|---|
| Model Link | View on HuggingFace |
| API Spec | View on HuggingFace |
| Github Link | No Github link provided |
| Paper Link | No paper link provided |
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Model overview
stable-diffusion-v-1-4-original is a latent text-to-image diffusion model developed by CompVis that can generate photo-realistic images from text prompts. It is an improved version of the Stable-Diffusion-v1-2 model, with additional fine-tuning on the "laion-aesthetics v2 5+" dataset and 10% dropping of the text-conditioning to improve classifier-free guidance sampling. This model is capable of generating a wide variety of images based on text descriptions, though it may struggle with more complex tasks involving compositionality or generating realistic human faces.
Model inputs and outputs
Inputs
- Text prompt: A natural language description of the desired image to generate.
Outputs
- Generated image: A photo-realistic image that matches the provided text prompt.
Capabilities
The stable-diffusion-v-1-4-original model is capable of generating a wide range of photo-realistic images from text prompts, including scenes, objects, and even some abstract concepts. For example, it can generate images of "a photo of an astronaut riding a horse on mars", "a vibrant oil painting of a hummingbird in a garden", or "a surreal landscape with floating islands and glowing mushrooms". However, the model may struggle with more complex tasks that require fine-grained control over the composition, such as rendering a "red cube on top of a blue sphere".
What can I use it for?
The stable-diffusion-v-1-4-original model is intended for research purposes only, and may have applications in areas such as safe deployment of AI systems, understanding model limitations and biases, generating artwork and design, and educational or creative tools. However, the model should not be used to intentionally create or disseminate images that are harmful, offensive, or propagate stereotypes.
Things to try
One interesting aspect of the stable-diffusion-v-1-4-original model is its ability to generate images with a wide range of artistic styles, from photorealistic to more abstract and surreal. You could try experimenting with different prompts to see the range of styles the model can produce, or explore how the model performs on tasks that require more complex compositional reasoning.
This summary was produced with help from an AI and may contain inaccuracies - check out the links to read the original source documents!
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