stable-diffusion-xl-refiner-0.9

Maintainer: stabilityai

Total Score

326

Last updated 4/28/2024

🗣️

PropertyValue
Run this modelRun on HuggingFace
API specView on HuggingFace
Github linkNo Github link provided
Paper linkNo paper link provided

Create account to get full access

or

If you already have an account, we'll log you in

Model overview

The stable-diffusion-xl-refiner-0.9 model is a diffusion-based text-to-image generative model developed by Stability AI. It is a Latent Diffusion Model that uses a pretrained text encoder, OpenCLIP-ViT/G. The model is not intended to be used as a pure text-to-image model, but rather as an image-to-image model to refine and denoise high-quality data. It is part of the SDXL model pipeline, which first uses a base model to generate latents and then applies a specialized high-resolution refiner model using SDEdit.

Model inputs and outputs

The stable-diffusion-xl-refiner-0.9 model takes an image as input and refines and denoises it based on the provided text prompt. It outputs the refined and denoised image.

Inputs

  • Image: An input image to be refined and denoised
  • Text Prompt: A text prompt describing the desired output image

Outputs

  • Refined and Denoised Image: The output image with improved quality and reduced noise

Capabilities

The stable-diffusion-xl-refiner-0.9 model is capable of refining and denoising high-quality images based on text prompts. It can be used to enhance the visual fidelity of images generated by other models or to improve existing images.

What can I use it for?

The stable-diffusion-xl-refiner-0.9 model can be used for research purposes, such as:

  • Generation of artworks and use in design and other artistic processes
  • Applications in educational or creative tools
  • Research on generative models
  • Safe deployment of models which have the potential to generate harmful content
  • Probing and understanding the limitations and biases of generative models

It should not be used for commercial purposes or to generate content that could be harmful or offensive.

Things to try

One interesting thing to try with the stable-diffusion-xl-refiner-0.9 model is using it in combination with the [object Object] model. The base model can be used to generate initial latents, which are then refined and denoised by the refiner model. This two-step pipeline can produce high-quality images while maintaining flexibility and control over the generation process.



This summary was produced with help from an AI and may contain inaccuracies - check out the links to read the original source documents!

Related Models

📊

stable-diffusion-xl-refiner-1.0

stabilityai

Total Score

1.5K

The stable-diffusion-xl-refiner-1.0 model is a diffusion-based text-to-image generative model developed by Stability AI. It is part of the SDXL model family, which consists of an ensemble of experts pipeline for latent diffusion. The base model is used to generate initial latents, which are then further processed by a specialized refinement model to produce the final high-quality image. The model can be used in two ways - either through a single-stage pipeline that uses the base and refiner models together, or a two-stage pipeline that first generates latents with the base model and then applies the refiner model. The two-stage approach is slightly slower but can produce even higher quality results. Similar models in the SDXL family include the sdxl-turbo and sdxl models, which offer different trade-offs in terms of speed, quality, and ease of use. Model Inputs and Outputs Inputs Text prompt**: A natural language description of the desired image. Outputs Image**: A high-quality generated image matching the provided text prompt. Capabilities The stable-diffusion-xl-refiner-1.0 model can generate photorealistic images from text prompts covering a wide range of subjects and styles. It excels at producing detailed, visually striking images that closely align with the provided description. What Can I Use It For? The stable-diffusion-xl-refiner-1.0 model is intended for both non-commercial and commercial usage. Possible applications include: Research on generative models**: Studying the model's capabilities, limitations, and biases can provide valuable insights for the field of AI-generated content. Creative and artistic processes**: The model can be used to generate unique and inspiring images for use in design, illustration, and other artistic endeavors. Educational tools**: The model could be integrated into educational applications to foster creativity and visual learning. For commercial use, please refer to the Stability AI membership page. Things to Try One interesting aspect of the stable-diffusion-xl-refiner-1.0 model is its ability to produce high-quality images through a two-stage process. Try experimenting with both the single-stage and two-stage pipelines to see how the results differ in terms of speed, quality, and other characteristics. You may find that the two-stage approach is better suited for certain types of prompts or use cases. Additionally, explore how the model handles more complex or abstract prompts, such as those involving multiple objects, scenes, or concepts. The model's performance on these types of prompts can provide insights into its understanding of language and compositional reasoning.

Read more

Updated Invalid Date

🗣️

stable-diffusion-xl-base-0.9

stabilityai

Total Score

1.4K

The stable-diffusion-xl-base-0.9 model is a text-to-image generative model developed by Stability AI. It is a Latent Diffusion Model that uses two fixed, pretrained text encoders (OpenCLIP-ViT/G and CLIP-ViT/L). The model consists of a two-step pipeline for latent diffusion - first generating latents of the desired output size, then refining them using a specialized high-resolution model and a technique called SDEdit (https://arxiv.org/abs/2108.01073). This model builds upon the capabilities of previous Stable Diffusion models, improving image quality and prompt following. Model inputs and outputs Inputs Prompt**: A text description of the desired image to generate. Outputs Image**: A 512x512 pixel image generated based on the input prompt. Capabilities The stable-diffusion-xl-base-0.9 model can generate a wide variety of images based on text prompts, from realistic scenes to fantastical creations. It performs significantly better than previous Stable Diffusion models in terms of image quality and prompt following, as demonstrated by user preference evaluations. The model can be particularly useful for tasks like artwork generation, creative design, and educational applications. What can I use it for? The stable-diffusion-xl-base-0.9 model is intended for research purposes, such as generation of artworks, applications in educational or creative tools, research on generative models, and probing the limitations and biases of the model. While the model is not suitable for generating factual or true representations of people or events, it can be a powerful tool for artistic expression and exploration. For commercial use, please refer to Stability AI's membership options. Things to try One interesting aspect of the stable-diffusion-xl-base-0.9 model is its ability to generate high-quality images using a two-step pipeline. Try experimenting with different combinations of the base model and refinement model to see how the results vary in terms of image quality, detail, and prompt following. You can also explore the model's capabilities in generating specific types of imagery, such as surreal or fantastical scenes, and see how it handles more complex prompts involving compositional elements.

Read more

Updated Invalid Date

📊

stable-diffusion-xl-base-1.0

stabilityai

Total Score

5.3K

The stable-diffusion-xl-base-1.0 model is a text-to-image generative AI model developed by Stability AI. It is a Latent Diffusion Model that uses two fixed, pretrained text encoders (OpenCLIP-ViT/G and CLIP-ViT/L). The model is an ensemble of experts pipeline, where the base model generates latents that are then further processed by a specialized refinement model. Alternatively, the base model can be used on its own to generate latents, which can then be processed using a high-resolution model and the SDEdit technique for image-to-image generation. Similar models include the stable-diffusion-xl-refiner-1.0 and stable-diffusion-xl-refiner-0.9 models, which serve as the refinement modules for the base stable-diffusion-xl-base-1.0 model. Model inputs and outputs Inputs Text prompt**: A natural language description of the desired image to generate. Outputs Generated image**: An image generated from the input text prompt. Capabilities The stable-diffusion-xl-base-1.0 model can generate a wide variety of images based on text prompts, ranging from photorealistic scenes to more abstract and stylized imagery. The model performs particularly well on tasks like generating artworks, fantasy scenes, and conceptual designs. However, it struggles with more complex tasks involving compositionality, such as rendering an image of a red cube on top of a blue sphere. What can I use it for? The stable-diffusion-xl-base-1.0 model is intended for research purposes, such as: Generation of artworks and use in design and other artistic processes. Applications in educational or creative tools. Research on generative models and their limitations and biases. Safe deployment of models with the potential to generate harmful content. For commercial use, Stability AI provides a membership program, as detailed on their website. Things to try One interesting aspect of the stable-diffusion-xl-base-1.0 model is its ability to generate high-quality images with relatively few inference steps. By using the specialized refinement model or the SDEdit technique, users can achieve impressive results with a more efficient inference process. Additionally, the model's performance can be further optimized by utilizing techniques like CPU offloading or torch.compile, as mentioned in the provided documentation.

Read more

Updated Invalid Date

🔎

sdxl-turbo

stabilityai

Total Score

2.1K

sdxl-turbo is a fast generative text-to-image model developed by Stability AI. It is a distilled version of the SDXL 1.0 Base model, trained using a novel technique called Adversarial Diffusion Distillation (ADD) to enable high-quality image synthesis in just 1-4 steps. This approach leverages a large-scale off-the-shelf image diffusion model as a teacher signal and combines it with an adversarial loss to ensure high fidelity even with fewer sampling steps. Model Inputs and Outputs sdxl-turbo is a text-to-image generative model. It takes a text prompt as input and generates a corresponding photorealistic image as output. The model is optimized for real-time synthesis, allowing for fast image generation from a text description. Inputs Text prompt describing the desired image Outputs Photorealistic image generated based on the input text prompt Capabilities sdxl-turbo is capable of generating high-quality, photorealistic images from text prompts in a single network evaluation. This makes it suitable for real-time, interactive applications where fast image synthesis is required. What Can I Use It For? With sdxl-turbo's fast and high-quality image generation capabilities, you can explore a variety of applications, such as interactive art tools, visual storytelling platforms, or even prototyping and visualization for product design. The model's real-time performance also makes it well-suited for use in live demos or AI-powered creative assistants. For commercial use, please refer to Stability AI's membership options. Things to Try One interesting aspect of sdxl-turbo is its ability to generate images with a high degree of fidelity using just 1-4 sampling steps. This makes it possible to experiment with rapid image synthesis, where the user can quickly generate and iterate on visual ideas. Try exploring different text prompts and observe how the model's output changes with the number of sampling steps.

Read more

Updated Invalid Date